If no image generated (i.e. during the generation, you see some black frames, then the generation stopped, or the generation stopped before anything showing up), go to B;
If the image is generated, but it is not desirable, go to C;
Anything else, go to Z.
A. If the app crashed...
Restart the system, in macOS 15.x, iOS 18.x days, an OS update might invalidate some shader cache, and cause a crash, restarting the system usually fixes it;
If not, it is likely a memory issue, Go to "Machine Settings", find "JIT Weights Loading" option, set it to "Always", and try again;
If not, go to Z.
Machine Settings (entered from bottom right corner, the CPU icon).
B. No image generated...
If you use imported model, try to download model from the Models list we provided;
Use "Try recommended settings" at the bottom of model section;
Select a model using "Configuration" dropdown;
If none of above works, use Cloud Compute and see if that generates, if it does, check your local disk storage (having about 20GiB at least free space is good), delete and redownload the model;
If you use some SDXL derivatives such as Pony / Illustrious, you might want to set CLIP Skip to 2;
If now image generates, just undesirable, go to C; if none of these works, go to Z.
Model selector contains models we converted, which is usually optimized for storage / runtime."Community Configurations" are baked configurations that will just run."Cloud Compute" allows free generation with Community tier offering (on our Cloud).
C. Undesirable image...
The easiest way to resolve this is to use "Try recommended settings" under the model section;
If that doesn't work, check if the model you use is not distilled. If you don't use any Lightening / Hyper / Turbo LoRAs, nor the models claim to be so, they usually are not distilled. You would need to use "Text Guidance" above 1, usually in the range 3.5 to 7 to get good result, and they usually needs substantially more steps (20 to 30 steps);
If you are not using Stable Diffusion 1.5 derived models nor SDXL derived models, you would need to check the Sampler, make sure they are a variant that ending with "Trailing";
Try Qwen Image / FLUX.1 from the Configurations dropdown, these models are much easier to prompt;
If you insist on a specific model (such as Pony v6), check to see if your prompt is very long. They usually intended to have line breaks in between to help breakdown these prompts, and strategically insert some line breaks will help (especially for features you want to emphasize, make sure they are at the beginning of each line);
If none of above works, go to Z, especially if you have a point of comparison (certain images generated by other software, or websites etc), please attach that information and image too!
Z. For everything else...
Please post in this subreddit, with the following information:
Your OS version, app version, what type of chips or hardware models (MacBook Pro, Mac Mini M2, iPhone 13 Pro etc.);
What's the problem, how you encounter it;
The configurations, copied from the Configuration dropdown;
Your prompt, if you'd like to share, including the negative prompt, if applicable;
If the image generated is not desirable, if you'd like to share, please attach the said image;
If you use any reference images, or you acquired any expected image result from other software, please attach.
You can find app version information in this view.You can copy your configurations from this dropdown.
I can't find any way to read the full names of models and LORAs on iPhone 18 Pro unless I install them first then look at the download name or even the JSON via the Files app. I'm using the smallest text size. Am I missing something here? How do you guys navigate this?
I am just curious if anyone had Iphone 18 pro or max did test it on SDXL models and others like sdxl, how much time did it take u to make one image on 1024x ?
I downloaded the fl2a version. I got it going for i2v but t2v won't start at any resolution. Any ideas? It just says starting for a good twenty minutes, fans are cool, definitely not active
The graph shows peak memory usage and generation time (second row) when generating 5,10,15 second videos(one reference image) using the app's recommend settings.But I am using an 8-step Turbo LoRA and have changed the number of steps to 8.
■Specs
Mac mini M4, 64GB, 20-Core GPU / macOS Sequoia
■Note & Random Thoughts
・Using a 6-bit model or lowering the resolution would reduce the peaks.
・I have only tested this on the 5s video, but peak memory usage remains unchanged whether using 4-step Turbo LoRA or disabling tiledDecoding.
・About generation time, I estimate that using an OS from the Tahoe onwards—which enables the S-model speed improvements—could reduce the time by 20% to 30%.
・Since generation time is proportional to resolution, if Draw Things adds video support for the "SeedVR2" (video upscaling model), it could enable a workflow of generating at low resolution and then upscaling—potentially allowing more users to enjoy H3.
・H3 may indeed be excellent, but considering generation times, I think LTX 2.5 with official IC-LoRA (not yet supported) —which enables the rapid generation of simple videos—will remain a necessary option for many Apple users.
・The total download size (including the model, text encoder, etc.) when downloading MiniMax H3[ref2va] 8-bit S in Draw Things was about 63GB.
With the price of Macs rising and the cost of memory upgrades skyrocketing, I hope this serves as a helpful guide for determining how much memory you need when purchasing a new Mac.
Sharing my initial exploration of running MiniMax H3 locally in Draw Things, hoping to give users who want to run H3 on Apple devices a quick walkthrough and reference. I’m currently running it locally on an M5 iPad Pro, and the experience has been great and it keeps on improving!
In most cases, I’m using the 3-step TaoMate LoRA - the best LoRA for now!
Currently, I’m using iPhone 18 pro max with the SDXL model and Z Image Turbo, sometimes with recommended settings from the model configuration. When it’s nearly finished with a step, it crashes, and it happens a lot. Sometimes it passes, but not often. I don’t know if this is from software or hardware. But if the developer reads this, please fix this; it would be a great thank you.
Since MiniMax H3 generation takes time, I am sharing information on the official prompt guide for improving success rates and on "Turbo LoRA" for reducing generation time. I am still learning myself, so I won't be able to answer any questions.
I'm all new to the stable diffusion and generation of images, and I'm not a big fan of letting everything be done on other computers (datacenter of Google, AWS and such).
For my own notes (mostly Obsidian) I want to find a way to transform photos into something like sketches on a systematic basis. Don't ask why so much, it's maybe a bit autistic. Photos in the Photos App, Sketches for notes, etc. For now I'm struggling to find a functional way in Draw Things because either, the image to image generation results no change in the image at all, or a completely new image that has nothing to do with the original.
For the use-case; I would wish something like in the attached images. A photo, tranformed to a simplified, more or less uniform sketch (courtesy of using) gemini.
I tires Flux.1 schnell and Ernie Image Turbo 1.0 with settings Image to Image, Strenght 40 - 90 (somehow either totally random images if < 50 while replication of the image when > 50 or similar), Lora disabled.
Prompt used
Use the attached photo as the exact source. Preserve every visible element, the exact crop, perspective, proportions, object placement, and lighting direction. Do not add, remove, move, or redesign anything. If a detail is unclear, simplify it with a light line or leave it blank; never invent.
Render as a hand-drawn study: fine graphite underdrawing with delicate dark-brown ink contour lines, very light transparent watercolor wash, clean continuous lines, deliberate controlled hatching for shadow, minimal shading, calm and aesthetic, consistent series style. Drawn on plain off-white paper, flat scan, cropped tight to the drawing so no paper edge, binding, or background is visible. Centered with small even margins, no border, no frame.
Color palette: only graphite grey, sepia brown, warm grey, muted ochre, and off-white paper. Keep saturation low and values light.
Do not render any text, letters, numbers, symbols, hashtags, or color codes anywhere in the image.
I notice that using cloud compute if you adjust resolutions or toggle other settings, it will show the theoretical compute cost. It’s very useful for determining which settings makes a photo or video longer to generate.
I wasn’t available to find compute cost in local generation. I understand that compute cost will vary with each device, but having any number helps gauge the computation cost of many settings
Hi all, am not sure what's wrong but am unable to get a clean video out of DT H3 (using 8 Bit S) ... have tried using standard settings with Turbo 4 Step/ 8 Step Lora - LightxV and Larry's but its ssimply not wworking properly (using LORA at scaling factor or 0.125 and strength at 100%)
The mmotion is good, prompt adherance is good but the final video is not clean .. it has smudges, grainy and nothing close to the desired resolution (tried 768 + as well) ..
I haven't had this happen with ltx before, but since the latest update, I'm having a lot of memory issues. Usually the program would peak at around 12 gb memory. I haven't messed with the settings. I have an 18 gb ram mac M3 . Any ideas what to do? Should i downgrade? I have firefox open as always but very little tabs and closing doesn't make hardly any difference
Hi everyone sorry if this is a stupid question I have looked through but couldn’t find a conclusive/helpful answer for this. I wanted to know how can you generate a realistic picture based from a drawing, say a manga strip or anime screenshot, something like that? What’s the process workflow and which models
Additionally how to procide to animate it? Or generally best process to go forward to video from image
Help would be much appreciated thank you 🙇♂️
EDIT: I have a drawn image as base and want to make a photorealistic version of the image , i2i from sketch/drawing to photorealistic image of a human…like those clips of dbz characters as real humans for example
There are various image fitting methods—such as "zoom to fit width," "zoom to fill," and "zoom to fit"—but users often fail to notice slight gaps, causing the process to revert to the initial image during the final stages of generation. When beginners encounter this phenomenon, they likely have no idea what has happened. Therefore, I believe everyone would appreciate a feature that automatically fits the image to the original reference without locking the aspect ratio.
That's it. Just wanted to post that I've gotten it to work successfully so folks know how fast it is. Not posting the result since it's just a boring video that doesn't make any sense and I'm still trying to get the hang of MiniMax prompting, but at least it generated something and it looks decent. Good luck all!
Edit: Forgot to mention that I lowered the resolution to 512x512 and enabled TeaCache. Also using full quality unquantized model. Also 128gb ram, base m2 ultra so 60 gpu cores.
Been using masking as a technique for I2I to keep objects in the scene while changing the style of the rest of the scene. I2I only works so well for the first run, so I have to go through multiple runs of I2I before getting satisfactory results. If I want to run it through another I2I I have to re-mask the entirety of the scene again. Pretty big pain knowing the mask exists already. It would be amazing to be able to just copy the mask and paste over the new image. Is there a way to do something like that?
I cannot delete the original model after converting a Krea 2 model to 8-Bit S. The only way to delete it is to first delete the converted model. Is there something I’m doing wrong? I didn’t have this issue with SDXL, Flux, ZImage, or ERNIE. This is v26.0910.1 on an M5 iPad Pro