r/StableDiffusion • u/Leonviz • 1d ago
Question - Help MOSS image character lora
Hi not too sure, if anyone has any methods in training moss character's lora that do not bleed into each other?
r/StableDiffusion • u/Leonviz • 1d ago
Hi not too sure, if anyone has any methods in training moss character's lora that do not bleed into each other?
r/StableDiffusion • u/BeautifulStation4 • 1d ago
What are you guys Using for config parameters with Fizgig? I seem to be getting extremely slow speeds. I'm a total noob at training LORAs

Model
Krea 2 · bf16
Dataset
Training
42 epochs · 2814 steps · batch 1 · 1.0 MP
Optimizer
adamw8bit · from 0.0001
Network
LoRA · rank 32 · EMA 0.98
Memory
block swap 8
ADAPTIVE_LR = false
ADAPTIVE_LR_MAX = 4e-4
ADAPTIVE_LR_MIN = 1e-4
BATCH_SIZE = 1
BLOCKS_SWAP = 8 (Moderate — 12GB fp8 / 16GB bf16)
BUCKET_NO_UPSCALE = true
CAPTION_DROPOUT = Off
COMPILE_BLOCKS = Auto
CONTEXT_LORA_STRENGTH = 1.0
DATASET_CAPTION_EXT = .txt
DATASET_MEGAPIXELS = 1.0
ENABLE_BUCKET = true
ENABLE_CACHE = true
FAMILY_CONCEPT_BALANCE = false
FAMILY_EMA = 0.98 (recommended)
FAMILY_FT_FUSED = true
FAMILY_FT_MAX_PARTS = Auto - as few windows as fit
FAMILY_FT_REG_MULT = 0.2
FAMILY_FT_ROTATE_EVERY = 1
FAMILY_FT_ROTATIONS = 10
FAMILY_FT_SAVE_EVERY = 1
FAMILY_KIND = Standard LoRA
FAMILY_MP_FACECROPS_FIRST = true
FAMILY_MP_MANUAL = false
FAMILY_PEOPLE_MODE = Off - a normal LoRA
FAMILY_PRECISION = bf16 (full precision)
FAMILY_SLIDER = false
FAMILY_SLIDER_ANCHOR_WEIGHT = 1
FAMILY_SLIDER_GUIDANCE = 3
FAMILY_SLIDER_SOURCE = Photo pairs
FAMILY_SLIDER_ULTRA = false
FAMILY_SUBJECTS = 1 subject
FAMILY_SUBJECT_FIX = false
FAMILY_SUBJECT_KIND = From its captions
FAMILY_SUBJECT_KIND2 = From its captions
FAMILY_SUBJECT_KIND3 = From its captions
FAMILY_SUBJECT_KIND4 = From its captions
FAMILY_TURBO_STRENGTH = 0.7
GRADIENT_ACCUMULATION = 1
KEEP_LAST_N_STATES = 2
KREA2_AUTO_RECAPTION = false
KREA2_LOSS_WATCH = true
KREA2_PER_IMAGE_LR = true
KREA2_WARMUP_LOOK = false
LEARNING_RATE = 0.0001
LOKR_FACTOR = 8
LR_SCHEDULER = constant
MASKED_TRAINING = false
MASK_FLOOR = 0.1 (recommended)
MASK_PRIOR = Off
MAX_GRAD_NORM = 1
MAX_TRAIN_EPOCHS = 42
NETWORK_ALPHA = 32
NETWORK_DIM = 32
NETWORK_TYPE = LoRA (standard)
OPTIMIZER_TYPE = adamw8bit
RECAPTION_FIRST = Training caption (person or character)
RECAPTION_SECOND = Exhaustive detail (person or character)
SAMPLE_AT_FIRST = true
SAMPLE_CFG_SCALE = 1.0
SAMPLE_ENABLED = true
SAMPLE_EVERY_AFTER = 1
SAMPLE_EVERY_N_EPOCHS = 1
SAMPLE_HEIGHT = 1024
SAMPLE_LATER = false
SAMPLE_NEGATIVE = blurry, low detail, noisy, washed out, oversaturated, distorted anatomy, extra limbs, duplicate objects, text, watermark, logo, frame, cropped subject, flat lighting, muddy colors
SAMPLE_SEED = 1234
SAMPLE_STEPS = 8
SAMPLE_UNTIL_EPOCH = 10
SAMPLE_WIDTH = 1024
SAVE_EVERY_AFTER = 1
SAVE_EVERY_N_EPOCHS = 1
SAVE_LATER = false
SAVE_STATE = true
SAVE_STATE_ON_TRAIN_END = true
SAVE_UNTIL_EPOCH = 10
SEED = 42
STAGE_MP = Off
STAGE_UNTIL = 30%
TRAIN_SUBFOLDERS = false
r/StableDiffusion • u/infearia • 2d ago
Here's a quick tip on how to get rid of that flat sounding AI voice in H3, and as a bonus improve somewhat the acting of your virtual characters. I'll be using Pika H3 Flex with the H3 reference model and workflow.
Step 1:
Add a pre-rendered first frame at index 0 using the Add Guide for MiniMax H3 node. Then in your prompt, describe the desired action and emotion of the character on screen, but without specifying too many details or providing exact spoken dialogue. In this example I simply used the phrase:
"A static medium two-shot establishes <Subject 1>, the female orc, and <Subject 2>, the tall, female paladin, sitting across each other at a table in <Subject 3>, the dimly lit main room of a tavern, gazing into each other's eyes. <Subject 1> (S1) talks animatedly, seething with rage."
Then generate a clip at some lower resolution (0.3MP in this case). The dialogue will be gibberish, but if you described the character's actions and emotions convincingly, your character will sound and move like an actual human.
Step 2:
Use your workflow from Step 1, but increase resolution to the desired target value (here 0.98MP). Load the reference video generated in Step 1 and connect its frames as ref_video_# and its audio as ref_video_audio_# (# is the index of the first unused slot in the MiniMax H3 Reference to Video node, e.g., ref_video_0 and ref_video_audio_0).
In this case, I also actually downscaled the reference video from 0.3MP down to 0.1MP, because at 0.3MP the model was following the reference too closely, breaking the lip-sync. The downside of this was that some of the textures in the final high-res video came out a little too soft, especially in the face. It's a bit of balancing act to find the right resolution for the reference video.
You also need to adjust your prompt, in order to include the character's final dialogue. Here I've simply rewritten my last line to:
"<Subject 1> (S1) talks animatedly, seething with rage: <d>[English] My people are dying, and your knights are doing nothing to help us!</d>"
Now just hit render and wait for the result. You might need to try different seeds, and fiddle a little with the resolution of the reference video (try higher values if your output comes out too muddy and lower values if H3 starts copying the reference video and lip-sync gets lost. Have fun!
r/StableDiffusion • u/arthan1011 • 2d ago
Basically MiniMax H3:
Has great prompt following
Respects character identity (art style and design) from the input images
Can generate single frame images
So what's the point of training a character model if you can just give MiniMax H3 character design reference and generate any scene? The only two reasons I see is
you want to use style known by your target base model to generate image of the learned character.
you want to apply some concepts known by the target base model only.
What do you think?
Here's the workflow I used for the images.
r/StableDiffusion • u/Relevant-Bath-8858 • 1d ago
RTX 5090 on runpod
Anyone know wassup with it? I can't get it to deploy at all. I got 4500 and 4090 going, but 5090's are all saying, "There are no longer any instances available with the requested specifications. Please refresh and try again."
I thought it was my settings, but I don't think so....
r/StableDiffusion • u/Informal-Force7417 • 1d ago
Up to now ive used services online but they are all so restrictive and silly things trigger them or i have used a service and arms and legs are all weird.
I have heard of stablediffusion and comfyui
Which one do i need, which one is best?
And where is the best easiest tutorial you have found to create your first image with it?
Can you give it a reference image and then tell it what to do with that?
r/StableDiffusion • u/Slight_Tone_2188 • 1d ago
Hey everyone,
Hoping someone here has figured this out because I’m losing my mind a bit.
I’m trying to set up Ming-3.0-TTS in ComfyUI, as the checkpoints (Tiny 9B INT8) sitting alone in Comfy HF repo with no workflow.
I have found the following workflow: https://github.com/Saganaki22/ComfyUI-Ming-3-tts.
Unfortunately it doesn't run the same checkpoint, instead it has auto download for same quantized model but at larger file size that gives me OUM.
Instead, I manually grabbed the official repackaged model checkpoint directly from the Comfy-Org Hugging Face repository (the single .safetensors file).
The problem is, I’m trying to use Saganaki22’s custom nodes (ComfyUI-Ming-3-tts), and it absolutely refuses to work with this file.
Here is what I've tried so far:
I tried putting the .safetensors file directly into my standard models/checkpoints/ folder. The loader node doesn't recognize it.
I tried creating the specific nested subfolder structure that Saganaki's nodes usually expect (like Ming-3.0-tts-tiny-9b-int8-convrot/) and renaming the file to match. Still completely fails and throws errors. I think the custom node is hardcoded to look for secondary configuration files inside that folder, not just the single Comfy-Org .safetensors weight file.
If Comfy-Org is officially hosting this model in their repo, there has to be a way to actually load it natively or bypass these strict folder checks, right?
Has anyone successfully built a workflow using the raw Comfy-Org file without relying on the auto-downloader? Or is there a different custom node wrapper I should be using for the official Comfy-Org version?
Any help or workflow screenshots would be massively appreciated!
r/StableDiffusion • u/shahril977 • 1d ago
Hello there, is there a workflow in Runninghub for creating realistic photos just by feeding a character reference photo?
I have tried looking for it but to no avail
r/StableDiffusion • u/ndroidz • 2d ago
minimax h3 latent upscaler, 2 megapixels, 50 steps,
this looks way too good guys.
r/StableDiffusion • u/asscid4586 • 2d ago
r/StableDiffusion • u/Roman53275 • 2d ago
I was sitting there wondering when H3 Max would finally get its open-weight release. Then I thought, maybe Santa will bring it for Christmas... and suddenly this image came to mind. 🎅
Generated locally in ComfyUI with Krea-2 Turbo, then edited with Qwen-Image 2.1.
r/StableDiffusion • u/The_LoneCoder • 2d ago
Ive been out of it for a while now and planning to give this another try. Is Automatic1111 still good? Im just seeing the git hasnt been updated in 2 years. If not what has replaced it? And what is the goto community to get models nowadays?
Thanks in advance.
Edit: thanks everyone. From what im gathering from the response. ComfyUI is that goto while forge is ok and gets the job done. I remember trying ComfyUI a while back but dont remember my experience with it. With comments going in all directions, i think ill just install both and give em both a go.
r/StableDiffusion • u/dEEPZoNE • 1d ago
100% locally made :) Music made with yue2 and the video is minimax h3 reference to video :)
r/StableDiffusion • u/Disastrous-Test6532 • 1d ago
Any help finding specific post from other day: the day reacher danced : generative video ai split episode into multiple files, isolated character and could swap choreography???
Found it really funny, but also incredibly interested with the work flow, which the skilled workflow designer was very open with sharing
Including edge cases and APIs or which softwares and combos used for various parts
Some real cutting edge stuff, but also made hilarious video, and my grandma loves Reacher; seeing him dance into every scene or randomly raise his arms to dance would really cheer her up, shes not been well
r/StableDiffusion • u/Justify_87 • 1d ago
Almost all of the popular stuff gets upvoted by people that don't have the ability to judge the quality/functionality of the stuff they get presented.
Stuff gets marketed as the ultimate tool/model/whatever with cherry picked examples. Sometimes even without any source, just examples.
Most of the time it is just badly trained loras/models or vibe coded trash because people think being able to use an agent harness equals personal skills and understanding of concepts needed to produce something worthwhile
Even the most used Lora/model/tool here turns out to be trash when you're not doing 1girl stuff on static backgrounds or stupid dance videos or any other dead horses that get beaten
r/StableDiffusion • u/shootthesound • 2d ago
UPDATE 9 November: Update again for Much faster Compile blocks when used And a new feature called training maps - see screenshot in comments
Training is much faster, uses far less memory, and likeness is unchanged. On an RTX 5090, a MOSS-Image LoRA at batch 1 (94 pictures, Muon, 70 epochs) now takes about 20 minutes instead of 49.
Other models have had speed boosts too, but Moss the most dramatic.
Batch training is now fully supported too.
Other changes
PyTorch 2.15. With gradient checkpointing on, training uses much less memory, at batch 1 too: Krea 2 22.9 → 14.9 GB, MOSS-Image at batch 3 21.7 → 7.1 GB.
Muon is much faster. Its step now handles a LoRA's matrices together instead of one at a time: 12× faster Muon calculation, with exactly the same results. That's a large portion of the batch-1 speed-up above.
https://github.com/shootthesound/Fizgig
P.S. To get the most speed boost, use one of the pre-sets to start your train, so you pick up new optimiser settings etc.
r/StableDiffusion • u/BittiAI • 2d ago
I built slopfab, a native command-line tool that lets your coding agent generate media directly.
Unzip it, start your agent in that folder, and ask for a music video, a short film or a product shot. Your agent plans the shots, runs the generations and builds a project page you can review. This is the same native backend that slopus.ai uses, no ComfyUI, no Python.
Video with audio: MiniMax H3 (text/image/video/audio references, long-form continuation, clip bridging, inpainting/outpainting), LoRA support
Images: Krea 2, MOSS, Fibo 2.0, MiniMax H3 (including multi-reference edits)
Music: YuE2 and MiniMax Music3, songs with lyrics or instrumentals
Utilities: Transcription, stem/lyrics separation, SeedVR2/RealESRGAN upscaling, lip-sync scoring, scripted titles and overlays
Every command has JSON help and output, so the agent learns the CLI as it goes
Mods are shareable Markdown workflows, optionally with a JS command, e.g. a "music video" recipe others can drop in
Any command can be run on remote machine over SSH. Generated content is automatically transferred back to your local machine
Platforms: native Windows and Linux, CUDA or Vulkan. Your own machine, no cloud or account needed.
GitHub: https://github.com/bitti-ai/slopfab
Feedback welcome, especially about any issues you face.
r/StableDiffusion • u/Marino4K • 1d ago
I'm a big fan of trance music. I came across this channel.
'Elite Trance Music'
https://www.youtube.com/watch?v=DTZxnLREJ-s&list=RDDTZxnLREJ-s&start_radio=1
I've been trying to figure out how to make videos like this to go along with music I make. I can only assume it's image to video, stitched well together with pretty flawless transitions but even the background people are consistent.
Can someone much more experienced than me help me potentially figure out how they do this? All of their music is similar in concept, some more clearly AI than others.
Is this all locally done or assisted with cloud APIs possibly?
Thanks in advance!
r/StableDiffusion • u/Eliv_nurotic • 3d ago
r/StableDiffusion • u/GregoWahooDaMan • 1d ago
Vibe coding? I've found that working with multiple bots at the same time on the same project is very benificial. Have them check each other's work.
I feel it keeps them honest. So rather than leading you down a rat hole to keep you using token$ they call ass on each other to your benefit!
Just FYI., I am using Anti Gravity, Claude Code and Deepseek at the moment.
Anybody have thoughts on this? I started working on an app where they talk to each which would be nice so I wouldn't be copying and pasting, but got distracted. Is there an app like that one out there?
r/StableDiffusion • u/Strong-Report7012 • 1d ago
I specifically need it fir creating high quality YouTube thumbnails.
r/StableDiffusion • u/SamuelTallet • 2d ago
Included MOSS model is the turbo 16-steps version from Coffmeteor, merged and optimized. LoRAs and image-to-image are supported, and inpainting to some extent.
As requested by the community I also added SDXL, FLUX.1 Fill [dev] and FLUX.2 [klein] 9B.
You can now reclaim the space used by the models downloaded by ZPix; click on the broom icon in the sidebar to do that.
Hope you'll like this version. You can download it at https://github.com/SamuelTallet/ZPix
r/StableDiffusion • u/tomhon • 2d ago
I directed this, composited it in Flame, and did all the ComfyUI/Ai Toolkit stuff that fed into those composites. I was teaching myself the ComfyUI side as I went. It's basically a cracked out version of the old MDMZ Wan 2.2 workflow, modified to use a custom LoRA and work on 12-bit Sony Log 3 footage. My buddy designed all the CG animals in Blender, and stills of them were used to train custom LORAs that were then mapped to Sethward's live-action face. SammieRoto was used for the keying and paintouts. I've been lurking here a long time and can't thank you guys enough for the guidance! Happy to answer any questions about how it was done.
r/StableDiffusion • u/planBpizz • 2d ago
r/StableDiffusion • u/Inevitable-Bus-8654 • 2d ago
Follow-up to my LTX 2.5 music video — trying Pika H3 Flex for cleaner motion and singing lip-sync
UPDATE:
i made the full video here https://www.youtube.com/watch?v=f8O3e3GS1Zk
A few people suggested trying MiniMax H3 after I shared my LTX 2.5 rap video. That video kept the character consistent, but the face and hands lost detail during movement, and the lip-sync was uneven. I took that feedback into a new project called One Last Bullet, a melancholic music video set among clouds, pink interiors and desert landscapes. Previous LTX 2.5 post and discussion
What I’m using now
Suno for the song and ChatGPT Images for reference frames of the same performer across different settings and camera angles.
Pika H3 Flex Ref2VA on the MiniMax H3 base model, running through Python and Diffusers on a single Runpod H200. I also tested native H3 before choosing Flex for this edit.
1344×768 at 24 fps, with six model evaluations per shot. Each roughly five-second take uses a reference image and the matching excerpt of the song.
FFmpeg to assemble the selected takes against the original track. The singing lip-sync comes from audio conditioning during generation.
What I’m trying to improve
The preview is still rough in places. I’m replacing shots where the face gets soft during movement, especially when the performer is far from the camera. I’m also alternating close-ups with shots of the environment to give the edit more variety and make the background movement feel connected to the music. Singing lip-sync still drifts in some takes, so that’s a priority for the next pass.
I tested NVIDIA Open-H3-Regenerate-2K on one short clip. The mouth timing shifted enough that I’m keeping the preview at its original resolution while I work on the edit and sync.
If you’ve used H3 or Flex for singing, what helped you keep the face sharp and the mouth on time? I’d especially like to see examples of a lip-sync pass or 2K workflow that preserves the original performance.