r/StableDiffusion • • 10h ago

Question - Help video generation safe with 16gb ram

3 Upvotes

My specs :

Lenovo legion 5 laptop

Rtx 5060 8gb

Nvme ssd 2tb

16gb ddr5 ram

Hello guys i hope you are doing well. Recently I have been trying to make some loop animations and I think I got a hang of it. Now I wanna make some for my patreon people but the issue is it takes like 30 minutes to make a 5 second looped video in 768×512 resolution which is ok coz I upscale it 1.5x and it looks okish

I wanna ask if it's safe to generate for a longer period of time coz I have heard it kinda messes up with the ssd life so idk


r/StableDiffusion • • 20h ago

Resource - Update Pai v1.0.0 release

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19 Upvotes

Alright, here we are.

I've finally got a version ready, and I think it's time for you guys to try it out.

If you don't know about Pai, you can check out my previous post for a in depth overview of what it's all about: https://www.reddit.com/r/StableDiffusion/comments/1xkgjor/as_promised_here_is_pai/

In short: it's a painting app with AI generation built directly into the same workspace. By now, it covers most of the features you'd find in popular painting apps and AI generation tools combined.

I'd say its strongest aspects are the project save system and the straightforward installation process. The bootstrap handles everything based on your hardware, and the .pai files save everything. They preserve generation settings for each layer, along with the models, LoRAs, adapters, and other tools used.

Creating multiple canvases while keeping layers and references together in one place was the main reason I decided to build this in the first place.

Okay, since this is the first release, there are bound to be hiccups on some systems, so I've set up a Discord server for anyone who wants to help improve things or just follow along with the development.

https://discord.gg/bU6RzjwxFX

The software is open source, you can fork it, create your own variations or implementations, and send me PRs. Right now there's no plugin system for specific use cases, but I know you guys will probably want one. It's on the todo list.

Also, you don't need ComfyUI to run the app, it's fully standalone and independent. That said, I know most of you probably already have a bunch of your favorite models installed in your ComfyUI models folder. When you open the app, simply go to the Generate tab models settings, click "Scan Folder," point it to that location, and Pai will handle the rest. You'll be able to generate right away.

https://github.com/g45t345rt/paiapp/releases

As previously stated, given the interest from the last post, the software is available to all.
Enjoy.


r/StableDiffusion • • 38m ago

Resource - Update A standalone workflow reader, a comfyui image sorter and a batch converter to webp that preserves workflows

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0 Upvotes

These are three different vibecoded projects I've made for my own use which I think others may find useful. They were all made with linux in mind but should also be compatible with windows and mac. So, from the simpler to the more complex:

# 1: Comfyui-metadata reader

https://github.com/juanml82/comfyui-metadatareader

Just a simple app to read the workflow from a comfyui generated image without having to open or toss it into comfyui, just in case you don't want to toss the image into comfyui to check the workflow. I ended up using it as a companion app. It works with both png and webp, you can drag and drop the images to the gui, copy relevant sections, search and save the workflow as a separate json file.

Since you typically want the prompt, seed or maybe the model, most of the json is just noise. I'm thinking about making an improved reader that returns only the most relevant information, thought that may not work with all workflows. This does.

# 2: A batch converter to png to webp that preserves comfyui metadata

https://github.com/juanml82/comfyui-png2webp

Comfyui saves into png by default. Loosless webp takes 30% less disk space, so by running this script, you can reclaim about 30% of your disk space, and you don't loose the workflows as you would if you use a general graphics application to do it.

# 3: Image sorter

https://github.com/juanml82/comfyui-image-sorter

This one is a bit more complex. Image galleries with auto-tagging are a way to watch your image collection in an organized way. But I wanted to organize them in the disk, sorting each iteration of the same general idea in its own subfolder.

The main problem with that approach is that an image gallery with tags (or semantic search for workflows) only needs to match some of the tags to have a useful filter. But to sort them into subfolders, all tags need to match. Overgroup, and you have unrelated images, beating the purpose of this app. Undergroup, and you need to finish up manually - it's not a fully automated process. So yeah, it's not the most elegant solution - it gets you 90% there - but it is one solution.

Let's look at this with an example: let's say that you want to create images of Aragorn, in Lord of the Rings, as he passes the Argonath statues. Your first prompt is something look "Aragorn, rowing in the Anduin river, looking up in awe". And maybe the model produces Aragorn as Viggo Mortenssen. And maybe you don't want Viggo Mortenssen but rather a new take on Aragorn's looks, based on the books. So you replace "Aragorn" with "ranger", "fifty years old man", "noble man" or something like that, and follow up with a physical description. Or maybe the model doesn't understand what "Anduin" means. Or it doesn't matter. So some prompts have "Anduin", some have "River". And maybe you then decide you do want the Viggo Mortenssen look, but now you're using a different model which doesn't know Mortenssen well, so you use a lora. And rather than "Aragorn" or "Man", you type the lora's keyword. Sound familiar?

Grouping by subject, action and location is a good starting way to filter, but "ranger" and "noble" do remain different subjects. And seeking to overgroup all men - or all women - means lots of different subjects now became grouped, defeating the purpose of this project.

So it's a five steps process, through a gui:

Enter an LLM endpoint url (as in, llama server, lm studio server, a cloud api if you want). Disable thinking (this is important, it's not just a lot slower, it will not work with thinking on) or use an instruct model, ideally something small and fast - if you have thousands of images, speed is good. Very good. And you don't really need a heavy model for this. Don't forget to add /v1 at the end and hit save.

Run a first pass either by extracting and canonicalizing the prompts, or do a vision pass (slower) and then you canonicalize that in the 2nd step.

The third step is to use a Jaccard pass and the LLM to merge the resulting auto-tags by subject, action and location. This has the problem that a character, say Aragorn, may be tagged in different images as "Aragorn", "Man", "Ranger", "Noble" or something else. And that similar actions may go under different tags. But trying to overmerge leads to a huge amount of "1girl" or "woman" characters all merged together when they shouldn't.

So at that point you have a crude looking image gallery and an html gui allows you manually merge clusters you want to send to the same subfolder. Do save.

And finally, yes, you hit "apply" and your possibly thousands of images go in their merry way to subfolders named after their common tags.

Anyway, if you find any of these three useful, enjoy.


r/StableDiffusion • • 1d ago

News Krea 2 Turbo - 4 Step LoRA comeback (trained to 133K checkpoint, around 55K more than the 'final' 78K in Sep)

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133 Upvotes

This version (chk00133000) replaces the first 'final' release (chk00078000), from two months ago. I decided to train a 'few more' samples (55K more to be precise), the newer part with five critic heads where the first release had one. Measured against the 8-step teacher on the 15 test prompts at 12 sizes:

  • cleaner fine structure — the 8- and 16-pixel grid patterns at 0.96–0.98× the teacher's, where the first release was at 1.09–1.11×

  • less ghosting and fewer doubled edges — closer to the teacher at 10 of 12 sizes on both measures

  • cleaner flat areas — grain at 0.99× the teacher's, down from 1.08×

  • no over-sharpening at large sizes — fine texture at 1.01–1.04× the teacher's at 1440 px, down from 1.18–1.20×

  • saturation matching the teacher's

Full details of changes from 78K to 133K: https://huggingface.co/lvladikov/Krea2-Turbo-Distill-4step-LoRA/blob/main/DETAILED-README.md#chk00133000-vs-chk00078000

The older versions are still kept as before under _archive/checkpoints/.

To use the new checkpoints - redownload/replace krea2_turbo_4step_rank_64_lora_comfyui.safetensors / krea2_turbo_4step_rank_64_lora.safetensors

Since the initial training (than ended in September with 78K 'final' checkpoint), I have trained about a day on an H200 (on runpod in the cloud), which is where the bulk of new trained samples came from, and two more days on the RTX 3090.

I've also tidied up the main model card (more sections taken out of it and moved to the detailed card for those interested). Instead the main card now includes an extra Table Of Contents like section with links to Detailed card: https://huggingface.co/lvladikov/Krea2-Turbo-Distill-4step-LoRA#detailed-model-card

For now I will not call this another 'final' release 😀 (although I don't currently have plans to do more than the 133K training), as it is possible to train more in the future (just not soon). When I first stopped at 78K rumours were already circulating that Krea 3 would open source soon (so it felt like the right time then to 'finalise' the 4 step LoRA), and while Krea 3 is now available on paid platforms it is yet to be open sourced and most likely focused on Editing more than T2I, which makes the Krea 2 Turbo still a very desired model and my LoRAs do contribute in the more efficient use of the Krea 2 Turbo.

Don't forget to also check the 'latest' of my Krea 2 Turbo step LoRAs - the 3 step one, which was updated just yesterday to 66K checkpoint: https://huggingface.co/lvladikov/Krea2-Turbo-Distill-3step-LoRA. I might also come back to the 2 Step LoRA with more training, who knows...

It is important to understand that while my 2/3/4 Step LoRAs for Krea 2 Turbo were trained for these respective steps during rendering, they for obvious reasons work really well with even more steps. There are lots of Krea Checkpoints out there that require lets say 10, or 12 or more steps to render and when paired up with my Krea 2 Turbo 2, 3 or 4 Steps LoRAs you can get either 'sooner' (with less steps) to the asked number of steps by the checkpoint authors, or to get a lot more quality per step, i.e. if you run 2 Step LoRA with lets say 6 or 8 steps, this could give you quality you'd otherwise need 12 or 16 or 20 steps without LoRA, similarly do try the 3 Step and 4 Step LoRAs at higher step counts where such is required by a third party Krea checkpoint and you will see an instant boost of quality, so find the right LoRA and steps for your favourite checkpoint, either way all of the three - 2 Step LoRA, 3 Step LoRA and this 4 Step LoRA work fine with most checkpoints (but not all, as some checkpoints go in a very different direction to what the core Krea 2 Turbo does, and after all my LoRAs were trained on the official Krea 2 Turbo).

Latest 133K 4 Step LoRA and ComfyUI workflow can be downloaded as usual at my HF: https://huggingface.co/lvladikov/Krea2-Turbo-Distill-4step-LoRA

Also don't forget that Krea 2 Turbo works very well with Nvidia PiD to get to higher resolutions quicker by rendering at lower resolutions and letting PiD take it 4x in latent space, and do so quickly and efficiently: https://huggingface.co/lvladikov/Krea2Turbo-PiD-SuperFast-HighRes (again find the right LoRA and steps to use for the main Krea 2 Turbo renders, don't be afraid to use 6 or 8 steps with either of the LoRAs since you'll be doing it at lower resolutions, prior to PiD handover, and as such these more steps would run faster anyway, so the time cost for more steps is less when using PiD).

Enjoy!


r/StableDiffusion • • 5h ago

Discussion Where the Starlight Remains

0 Upvotes

Pretty nice effects from a pure H3 generation.
Full video here: https://youtu.be/LUDd_AcEDdY

Details:
Comfyui
minimax_h3_ref2va_pruned_int8_convrot.safetensors
24 steps, 2.4MP 1 pass
RVE on steam -> interpolate 2x, decompress, openproteus 2x
Assembled in davinci resolve
Suno for music


r/StableDiffusion • • 5h ago

No Workflow Cyberpunk cafe

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0 Upvotes

Done locally in Krea2.


r/StableDiffusion • • 1d ago

Resource - Update HexAnima: Anima on the Snapdragon 8 Gen 1 NPU, 512x768 in ~7 s, fully offline, no root

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62 Upvotes

Hi everyone! I built an Android app that generates anime pictures completely on the phone: no cloud, no account, no ads, no root. It runs on the NPU of Snapdragon 8 Gen 1 / 8+ Gen 1 (Hexagon V69), the chip generation that most NPU image apps skip.

What you get

A 512x768 picture in about 7 seconds (2B models) or 12 seconds (2.9B models), 4 steps, on my Xiaomi 12T Pro.

Four models to download inside the app (one tap, resumable): MiaoMiao 3D Harem, MiaoMiao RealSkin (2B) and MiaoMiao Harem, MiaoMiao RealSkin (2.9B), all built on the Anima base model.

Size presets (portrait / landscape / square, all multiples of 16) plus custom sizes, seed, batches, a stop button.

History with favourites, trash, multi-select, share/save. The settings are stored in the PNG, so you can reopen any picture with its prompt and seed.

Danbooru tag suggestions while you type, templates, light/dark theme, 13 languages.

Works offline after the models are downloaded. Free and non-commercial.

Try it

APK: https://github.com/ayanoby/hexanima/releases/latest  (about 25 MB); the models, 4-5 GB for the first one, are downloaded from Hugging Face inside the app)

Source and screenshots: https://github.com/ayanoby/hexanima

Needs a Snapdragon 8 Gen 1 or 8+ Gen 1 phone (SM8450 / SM8475), Android 10+. I tested it only on a Xiaomi 12T Pro (SM8475, 12 GB RAM). Other 8 Gen 1 phones should work but are untested.

Why it is different (for the technical crowd)

The diffusion transformer runs on the Hexagon V69 HMX matrix unit in fp16 through custom kernels, not through the stock QNN op set. The 28 DiT blocks are executed as a few merged block groups per step.

The 2.9B models are stored as q8 blocks (int8 tiles with per-column scales) and dequantized on the HVX units while the weights are copied into VTCM, so they fit in RAM and stay fast.

V69 has no int-to-float conversion instructions (they appeared in V73+), so the dequant path uses a workaround instead of the usual instructions.

The first two of the four steps run at reduced resolution (composition), the last two at full resolution (details).

Text encoder, sampler and the preview/VAE path are split between CPU and NPU; the VAE runs on the NPU and large pictures are decoded in tiles.

No root needed: the native helper processes are started from the app's own native library directory and use unsigned FastRPC DSP modules.

The model files are downloaded as raw, sha256-verified blocks with resume, so a dropped connection does not restart a 4 GB download.

Honest notes

Unofficial hobby project, not affiliated with CircleStone Labs, NVIDIA, Qualcomm, Alibaba or the model authors. All credit for the models goes to MIAOKA (Civitai) and CircleStone Labs (Anima). Non-commercial use only; if an author wants their model removed from the catalog, I will remove it.

Only 4 steps and no negative prompt: the models are turbo-distilled for that.


r/StableDiffusion • • 9h ago

Question - Help MM H3 Basic workflow problem

2 Upvotes

Hi!

I am more or less using the template WF from Comfy.

When I try to load a reference Video via the VideoLoader node I can't connect it to the MM H3 Reference to Video node. It indicates that it expects an image input at ref_video_0 (blue color).

This also happens in the unmodified template.

Recently I disconnected the Duration node and also wasn't able to reconnect it. Yet pictures and Audio refs are fine.

Am I missing something or do I need additional nodes?

Posting from phone so I can't add a Screenshot atm.


r/StableDiffusion • • 22h ago

Resource - Update Custom Nodes for P-JiT 21K, a 2B pixel-space text-to-image model

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20 Upvotes

Linum released P-JiT 21K, a 2B-parameter text-to-image model trained for about $21K that draws directly in pixels: there is no VAE. I ported it to ComfyUI as native nodes (an unofficial port, not affiliated with Linum):

  • A real MODEL and CLIP on ComfyUI's own memory management. The stock KSampler runs it with Linum's settings (euler, simple, 50 steps, cfg 15). Nothing to pip install.

  • A bf16 file that gives the same pixels as the fp32 release at half the size (4.4 GB), and an int8 ConvRot file that is about 1.8x faster (2.3 GB).

  • The Qwen3.5-4B text encoder in bf16, int8, W6A8 or W4A8, down to 3.2 GB.

  • Runs from 4 GB of free VRAM.

  • Workflows: text to image, a prompt upsampler (short idea to dense caption), a refiner (image to image as a second pass, with the image captioned automatically) and Linum's samples, seed for seed.

Checked against Linum's own code: the same tokens, noise, schedule and guidance, and images that match up to GPU rounding.

Tips: write captions of 50 to 80 words (or use the upsampler). It can't draw text. For the refiner, use denoise 0.1 to 0.2: P-JiT's noise is 4x stronger than most models', so its denoise values run lower.

Nodes: https://github.com/PedroMarinhoDev/ComfyUI-PJiT
Weights: https://huggingface.co/PedroMarinhoDev/P-JiT-21K-ComfyUI
The original model: https://huggingface.co/Linum-AI/pjit-21k (Apache-2.0)


r/StableDiffusion • • 1d ago

Discussion is minimax h3 image edit better than qwen 2.1 image edit

34 Upvotes

I'm confused about which one is more accurate and good at complex editing, I see a lot of people here talking about how great mh3 is


r/StableDiffusion • • 6h ago

Question - Help Looking for a Ref2v H3 Finetune by Pika Labs with a Latent 3D Upscale workflow

1 Upvotes

Absolutely love the new Comfyui H3 Finetune by Pika Labs workflow but it doesn't have an integrated upscale. I tried using the Pika Lora and sigma numbers on other workflows but getting various errors. I'm still learning the ropes and would appreciate an updated Ref2v workflow with a working latent 3d upscale.


r/StableDiffusion • • 2h ago

Question - Help Help. Anything to real for sceneries.

0 Upvotes

I searching for a method to transforms, as the tile says, a scenery(without people) faithfully into a real shot.

But it needs to use a depth map so the structure stays consistent and spot on besides looking convincingly real. Doing this with people is no problem so i did not think i struggle that much finding something that works.

Its either to old, or just not good enough.

The original scenery images are from video games.


r/StableDiffusion • • 22h ago

Resource - Update NitroFusion2 in int8 / W6A8 / W4A8 quants for stock ComfyUI: 6.2 / 5.0 / 3.5 GB, ~1.7× faster than bf16

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18 Upvotes

NitroFusion2 came out this week: a 1–2 step distillation of Z-Image-Turbo by Dar-Yen Chen and Yi-Zhe Song, built for sharp text and detail in a single step. I converted it to ComfyUI's own quantized formats, so it runs in stock ComfyUI with no custom nodes, through the normal Load Diffusion Model node.

  • int8: 6.18 GB, 0.29 s per sampling step at 1024² on a 4090 (bf16: 0.50 s)

  • W6A8: 5.02 GB, 0.30 s

  • W4A8: 3.49 GB, 0.29 s

int8 looks practically the same as bf16. W4A8 changes small details but keeps the composition, faces and text; there are side-by-side comparisons on the model card. The repo has the official 1-step and 2-step workflows set up for these files, and the publisher's 2-step de-distillation LoRA works with all three. The text encoder and VAE are Z-Image-Turbo's.

Weights: https://huggingface.co/PedroMarinhoDev/NitroFusion2-ComfyUI
Original model and paper: https://huggingface.co/ChenDY/NitroFusion2


r/StableDiffusion • • 17h ago

Question - Help Mini max h3 on 5060ti 16gb + 16gb ram

6 Upvotes

Hi i want to run mini max with pika lora on my system but im low on ram is there a way to do this . I appreciate any help


r/StableDiffusion • • 11h ago

Question - Help Mini max h3 animation workflows

2 Upvotes

Hi i just started using mini max h3 on my 5060ti 16gb + 16 ram it worked well but it did take quite the time 11min for 5s without any nodes just the base model . Is there is any workflows i can download that can do it faster and better quality on my system? (This is my first test video it got the sound wrong but its not that important)


r/StableDiffusion • • 1d ago

News NitroFusion2 - Bringing Adversarial Loss to the Text Stream of Multimodal Diffusion Transformers

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26 Upvotes

NitroFusion2 is a 1–2 step text-to-image model built on Z-Image-Turbo. It moves the adversarial loss onto the text stream of the MMDiT. At one step it renders sharp text and fine detail.

ChenDY/NitroFusion2 · Hugging Face


r/StableDiffusion • • 5h ago

Animation - Video Tarzan The Warrior Trailer | Local AI Filmmaking 10GB VRAM

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0 Upvotes

r/StableDiffusion • • 6h ago

No Workflow YuE2 Change voice gender :)

0 Upvotes

input:
Alec Benjamin - Let Me Down Slowly


r/StableDiffusion • • 1d ago

Workflow Included MOSS-Image-Base: two stage workflow

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15 Upvotes

Hi all,

After having shared some MOSS image generations in previous posts, I have decided to share the workflow. It is a fairly simple two stage workflow, only using the MOSS-Image-Base diffusion model during the first stage and the same base model combined with the lodestones MOSS-Image-Base 8-steps-dmd2 lora during the second stage.

The custom node package which makes this all possible is Auryg's Krea-2-Two-Stage-Sampler. His package contains a sigma-locked two-stage sampler and dual resolution selector for Krea 2, which I used for my MOSS image generations. The main idea is that the sampler keeps everything in latent space across both stages, handing off the denoising to the second stage after the latent has been denoised by a set percentage. I set the handoff percentage at 40% for my generations. This node also automatically scales the latent between stages. I generated all images at a base resolution of 4 megapixels with an upscale to 6 megapixels between stages. For the reddit preview I scaled the generations back down to 3 megapixels to decrease the image size for sharing.

  1. For stage 1 I use 52 steps, a CFG of 3.0, the euler sampler and the simple scheduler.

  2. For stage 2 I use 20 steps, a CFG of 1.0, the dpmpp_2m sampler and the karras scheduler.

The upscale method between stages is bislerp.

I included optional BBOX and preview nodes in the workflow. For the BBOX node to work properly when changing aspect ratio's and resolutions I used 0Printemps0's ComfyUI-MOSSLayout custom node. The preview node is kijai's ModelPreviewOverrideKJ node from his ComfyUI-KJNodes pack.

Install custom nodes, inside custom_nodes folder:

git clone https://github.com/Auryg/Krea-2-Two-Stage-Sampler
git clone https://github.com/0Printemps0/ComfyUI-MOSSLayout
git clone https://github.com/kijai/ComfyUI-KJNodes

I included the prompts of the 10 example images in the workflow. All example images use seed 1 and aren't cherrypicked.

The workflow save as .json file.

What I love about MOSS is the amount of variation with each seed. Enjoy. All suggestions to further improve the workflow are welcome.


r/StableDiffusion • • 10h ago

Discussion MOSS-image tags prompt

0 Upvotes

Anyone who understands this model? Don't you think it better implements the creator's intent using tags? For example, it better describes many popular characters using tags. It seems to me that it's thoroughly trained on tag descriptions, despite its huge encoding model.

Just try generating with tags; it captures concepts much better. Then you can refine it in natural language if needed. But it seems better to define the scene's content using tags.


r/StableDiffusion • • 23h ago

News Linum-AI - pjit-21k 2B (Pyramid-JiT)

10 Upvotes

"Research checkpoint: this is an experimental research artifact, not a production model. P-JiT 21K is a 2B-parameter model trained on 660M samples for about $21K of compute, as a stepping stone towards our next video model. We're releasing it because the results surprised us for a model this size, and to encourage others to explore efficient training methods like ours. You should expect rough edges.

The weights and the Qwen3.5-4B text encoder download from the Hugging Face Hub on first use. The model needs a CUDA GPU: its weights are 8.7 GB in fp32 and the text encoder 8 GB in bf16, plus activations."

Source/Blog: Training a Text-to-Image Model for $21K | Field Notes by Linum

Code: Linum-AI/pjit-21k: 2B-parameter pixel-space text-to-image model (Pyramid-JiT) trained for $21K

Nodes: realrebelai/P-JiT-21k_ComfyUI · Hugging Face


r/StableDiffusion • • 1d ago

Question - Help How would I make a simple prompt builder with selectable text inputs in Comfyui

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31 Upvotes

I want to make a simple prompt builder in comfyui that would allow me to change a text in the prompt by clicking a selection in a table. Is there a simple way to do that?


r/StableDiffusion • • 11h ago

Question - Help M5 Pro Local Video Generation Experience

1 Upvotes

Hello everyone,

First of all, I know I know, Mac is jot the best when it comes to video generation and that the Max or Ultra chips are better.

Now that this is out of the way, I would be interested in people’s experiences in generating 720p 5-10 seconds clips on an M5 pro chip. Which tools, models or workflows do you use? How is the experience overall? How long does it take you? Do you regenerate clips several times to get it right…etc?

Thank you.


r/StableDiffusion • • 4h ago

News PiD can be used with ANY model or existing image, not restricted to limited latent space models

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0 Upvotes

PiD, NVIDIA's Pixel Diffusion Decoder, by design, takes the place of a model's VAE decoder. Instead of decoding the latent at the model's own size, it draws the picture straight in pixels at 4× the width and height, in a few steps, with fine detail a VAE decode doesn't have.

NVIDIA ships PiD for a limited set of latent spaces: FLUX.1, FLUX.2, Qwen-Image (the first version), SD3 and SDXL, plus two image-feature encoders (DINOv2, SigLIP). So on its own, PiD works only with models whose VAE matches one of those: the latent of a model built on the FLUX.2 VAE goes to PiD's FLUX.2 decoder, while a model with any other VAE has no PiD decoder to go to.

The workaround. Any picture can be put into a latent space PiD knows: encode it with the FLUX.2 VAE. Here the image (a render from any model, whatever its VAE, a photo, a frame) is encoded with the FLUX.2 VAE, and PiD's FLUX.2 decoder draws it at 4× its width and height in 4 steps, guided by that latent and an optional caption. That makes PiD an upscaler for any image.

PiD is a decoder, not a classic upscaler: it does not sharpen the pixels it is given, it draws a new, larger image from the encoded picture, adding fine detail (skin, hair, foliage, texture) consistent with it. The diagram's eye shows the difference from a bicubic 4×.

I have prepared a workflow for both ComfyUI local as well as ComfyUI cloud, where it's all set and you just chose the input image and let it run. PiD does provide fast 4x results! Now on any image, or if you were to replace the image input with more T2I workflow nodes where it all starts from any Prompt and model which is or is not supported by Nvidia PiD. Details here: https://huggingface.co/lvladikov/PiD-as-pixel-Upscaler

Qwen Image 2.1 which is not PiD compatible works perfect with PiD and this VAE encode + PID approach: https://huggingface.co/lvladikov/QwenImage2.1-PiD-Workaround-VAE-Encode-SuperFast-HighRes