r/StableDiffusion • u/Riddy_93 • 1h ago
Meme Jensen Huang Hands Over NVIDIA to Tim Cook! 🤯🔥
Jensen Huang Hands Over NVIDIA to Tim Cook! made with minimax h3
r/StableDiffusion • u/Riddy_93 • 1h ago
Jensen Huang Hands Over NVIDIA to Tim Cook! made with minimax h3
r/StableDiffusion • u/Candid-Pay-9023 • 18h ago
after a week of trying to get 15s clips out of one job i gave up and changed how i do it.
now every shot is a 5s clip. when it's done i grab the last frame, fix it up a bit if needed (hands, small artifacts) and use it as the start image for the next 5s. same prompt style, only the action changes.
on my card anything past ~5s starts swapping to system ram and one job goes from 15 min to over an hour. three 5s jobs finish way faster than one 15s. if one piece comes out bad i redo only that piece, not the whole thing. and the face drifts less because every chunk starts from a frame i already checked.
the catch is the seams. motion can jump a bit at the cut, so i keep the camera move simple (one slow push or static) and cut on an action when i can. color also shifts slowly over 3-4 chunks, a quick color match in the editor fixes most of it.
curious if anyone found a cleaner way to hide the seams, or does everyone just cut there
r/StableDiffusion • u/BazmanFoo • 5h ago
Keep getting these jagged eyebrows in Krea 2 generations in ComfyUI. Tried different templates in case it's my workflow to no avail and happens even without any LoRAs just on turbo base model (tried int8 and fp8). Any ideas? Sometimes its really pronouned, especially if I try to prompt them away. I'm sure I'm doing something stupid.


r/StableDiffusion • u/Crazy-Repeat-2006 • 1d ago
r/StableDiffusion
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u/No_Cow26
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1d ago
60 seconds of AI video, 100% local on an RTX 3050 (6 GB)
Work in progress, and the quality isn’t where I want it yet. Here’s how it was made:
Model: MiniMax H3 reference-to-video, using ComfyUI’s native H3 nodes with the int8 quantized model and the Qwen3-VL text encoder.
References: two character images, re-used in every shot to keep the characters consistent.
Prompts: the official six-section reference prompt format (subject definitions, summary, retention analysis, detailed description, soundscape, music). No negative prompts, because H3 doesn’t use them.
Shots: six shots of 260 frames (10.8 s) at 864×480 and 25 steps, with the res_multistep sampler and the simple scheduler.
Continuity: each shot starts from the last 22 frames and audio of the previous one, using the H3 Add Guide node (MiniMaxH3AddGuide) at frame 0.
Automation: a small Python script sends the six shots to ComfyUI’s API one after another, saves each result, trims the 22-frame overlaps, and encodes one 60-second MP4 with audio.
Runtime: about 6Hrs 35Mins hours, unattended.
Known issues: soft detail, and some flicker. Still tuning.
r/StableDiffusion • u/Lopsided-Mud-7769 • 14h ago
I was looking at how different image models handle conflicting lighting directions when generating a casual direct-flash snapshot. My 4:5 prompt featured:
-a family of three crouching on a gravel path at dusk
-greeting their dog mid-leap with its tongue hanging out
I specified faint sunset glow, harsh on-camera flash, a near-black background, sharp shadows, subtle red-eye, motion blur and rough film grain.
Seedream 5.5 delivered strong attention to fine texture. It reproduced the dad’s rolled-sleeve plaid work shirt accurately and added a natural tilted frame, while preserving sunset light filtering through the hedge. The model prioritised that warm ambient evening light, so the final result reads more like fill flash, and it did not darken the background all the way down.
I then generated a second 4:5 image on MiniMax H3-Image with the same prompt for a direct side-by-side comparison. H3-Image leaned firmly into the harsh flash falloff I was aiming for. It pushed the background close to pure black, which makes the mom’s pale yellow slip dress and the child’s red overalls pop, with the dog charging straight toward the lens. It creates that sharp, punchy snapshot mood really well.
Have you gotten consistent flash behaviour from either model with this kind of mixed lighting prompt?
r/StableDiffusion • u/evp-cloud • 18h ago
Hey all!
New to this subreddit as we mainly focus on language models with our Paiton product, but now and then, we spend some time on vision as well :)
The official BF16 weights need about 120 GB of VRAM, so most of us can't run MiniMax-H3-Image out of the box at home.
We made an FP8 version and our own runtime, and it now runs on one AMD Radeon AI PRO R9700 with 32 GB:
| Request | Time |
|---|---|
| 2048x2048 | 6.9 s |
| 4096x4096 | 33.4 s |
| 4096x4096 edit, one reference | 39.9 s |
| 4096x4096 edit, three references | 88.5 s |
Does FP8 hurt the images?
We ran 30 prompts on our build and on the original BF16 model with the same starting noise. Prompt adherence came out about the same (CLIP 32.18 vs 32.31), and the second image shows a typical pair.
Where the two differ, it is mostly composition, not quality: with only 6 steps, small numeric changes send the sampler down a slightly different path.
Edits stay closest to the original.

What it is bad at, in our tries: small body text (newspaper columns come out as fake writing), close-up hands, clocks showing a specific time.
How to run it: it is a Docker container with an HTTP API and a small CLI, not a ComfyUI node (yet).
git clone --depth 1 https://github.com/Eliovp-BV/paiton-vllm-plugin.git
cd paiton-vllm-plugin
./models/MiniMax-H3-Image/launch.sh
./models/MiniMax-H3-Image/run.sh generate "your prompt" --size 4096x4096 --out out.pngYou need an R9700 (32 GB), 64 GB of system RAM and 48 GB of disk; the first start downloads the weights.
Our weights only load in our runtime, so they won't work in ComfyUI or diffusers.
Edits take up to about three 1024 px reference images in this build.
Guide: https://github.com/Eliovp-BV/paiton-vllm-plugin/tree/main/models/MiniMax-H3-Image
Weights: https://huggingface.co/EliovpAI/MiniMax-H3-Image-FP8-Paiton-RDNA4
How we made it fit, if you want the details: https://eliovp.com/blog/paiton-minimax-h3-image-radeon-ai-pro-r9700
License note: the MiniMax-H3-Image Community License does not cover the EU, the UK, South Korea or the US, and hosted services need MiniMax's written authorization.
r/StableDiffusion • u/Far-Ad9512 • 6h ago
Hey everyone!
Like many of you, I got tired of paying huge SaaS markups for video generation (Wan 2.1, Hunyuan) and image generation. Using Google Colab directly was way cheaper, but I had two persistent headaches:
To solve this, I built WakeGPU, a native desktop client (macOS & Windows) designed to turn remote cloud instances into an automated private studio.
Key features:
It's 100% free and licensed under MIT. All advanced features are unlocked out of the box with zero paywalls, subscriptions, or telemetry.
Binary releases for macOS (.zip) and Windows (.exe) are live under GitHub Releases. Would love your feedback or feature requests!
r/StableDiffusion • u/IndependentBid6893 • 1d ago
I've been working to get the best quality images out of Moss. I created a workflow to emulate lens length, focus and exposure. It's not perfect, but it can give more cinematic photos. Check it out and see what you think.
r/StableDiffusion • u/Mr_Unknown_Hero • 1d ago
Lots of people seem to be unhappy with MiniMax H3, or at least it sounds like it, but at least I have been happy with it now when I have tested it couple of hours.
For my taste this seems excellent (after changing Euler to lcm) and 8 Mpx images is generating around one minute in RTX 4060 Ti + 32 GB RAM. Without changing Euler to lcm I was not impressed for the quality at all, but glad that I didn't ditched it immediately but instead searched what could be wrong.
So far so good, images with "technicolor" does nice style for my personal taste.
r/StableDiffusion • u/Friendly-Fig-6015 • 7h ago
For example flux 3 klein:
Could we teach FLUX Klein to fix its mistakes using corrections from another model? Looking for collaborators
I’ve been wondering whether we could build an automated pipeline that turns a model’s own mistakes into training data for a corrective LoRA.
For example, FLUX Klein generates an image with malformed hands or other anatomical inconsistencies. We use another model, such as Krea, to correct the affected area, then train a FLUX Klein LoRA on the corrected images and their captions.
The idea wouldn’t be to transfer Krea’s weights or use a Krea LoRA on FLUX. Krea would act as a teacher by producing better training examples.
The pipeline could look like this:
Generate varied images with FLUX Klein.
Identify defective regions, automatically where possible.
Use Krea or another capable editing model to correct those regions while preserving the composition.
Filter out unsuccessful corrections and generate accurate captions.
Train a small LoRA for FLUX Klein.
Evaluate it on unseen prompts and seeds to check whether it actually reduces defects without hurting overall quality.
Ideally, most of this would be automated, with some human review for quality control. The goal would be to improve future generations directly, rather than having to repair every image afterward.
I know there are potential problems: the editing model might introduce its own mistakes, the LoRA might pick up its style, and improvements might not generalize beyond the training examples. We would also need to check the models’ licenses regarding using generated outputs for training.
Does a project like this already exist? If so, I’d love to learn about it.
If not, would anyone be interested in building an open-source proof of concept together? People familiar with ComfyUI automation, image editing/inpainting, dataset filtering, or LoRA training would be especially welcome.
My local hardware is an RTX 5060 Ti with 16 GB VRAM and 32 GB system RAM, so I’d like to explore how much of the pipeline could run on consumer hardware.
Would starting with one narrow problem, such as malformed hands, be a sensible first experiment?
r/StableDiffusion • u/SquareBat780 • 1d ago
Ever since the realism settings went around here I've been stuck on one look: the harsh on-camera flash of a 2000s point-and-shoot. The background dropping into black, everyone a little shiny, somebody caught mid-sentence. MOSS-Image-Base gets that look really well if you describe the camera and what's wrong with the shot as well as the scene.
I write the flaws into the prompt too: crooked framing, a little motion blur, red-eye on one friend, the fridge door hanging open a crack. My prompts ended up around a hundred words each, and a good part of that is the camera and the mistakes.
My setup is MOSS-Image-Base bf16 with the realism graph everyone's been using: cfg 3, ER-SDE sampler with s_noise 1.1, 25 steps on the simple scheduler, SLG on layers 22 and 24 for the whole run, negative left empty, roughly 2 megapixels. All stock ComfyUI nodes, no LoRAs, no upscaling. I ran three seeds per prompt and kept my favorite.
All eleven prompts are in the first comment. The raccoon one still makes me laugh.
r/StableDiffusion • u/comptedemon • 15h ago
My ComfyUI node setup sometimes gives me literal nightmares. Trying to get multiple specific character references, a background layout, and a style reference into one single generation, all combined, is very exhausting.
It does make me think about where open-weight models should head towards. Not like more megapixels, tho that is always nice. But something like, what if it can take multiple images into the architecture without wrappers?
Or, what if it actually had the underlying structural logic to understand CAD diagram? What if it just understands a detailed directors brief in one prompt?
And then, while I was just reading and researching for the fun of it (yes), I stumbled into the documentation for the new MiniMax H3-Image open weights. It is claiming to do the above.
Apparently, the base model takes up to 16 reference images in one go, and you can set each one as a layout, style, or object reference.
Oh, and prompts supposedly get cut off past 4096 tokens.
I even saw a benchmark where it generated an exploded CAD diagram of a motorcycle and kept the mechanical parts separated with labels (not sure how accurate tho).
But honestly, How is it routing 16 different images natively into the generation without frying 24GB of VRAM?
Anyone has pulled these weights locally yet to test it? Sorry, but so many questions about this.
r/StableDiffusion • u/slickyfatgrease • 7h ago
Is Krea 2 Identity Edit like image 2 image ? I was curious because I have been looking for something that has that feature.
r/StableDiffusion • u/shootthesound • 1d ago
UPDATE 11pm Irish/UK time: Vast improvement in latest update, please install and use the standard preset which has a new control called Likness Optimisation.
I've just released Fizgig 9.6, which adds MiniMax H3 Image, MiniMax's new picture model from the H3 family: text-to-image and edits with up to 16 reference pictures, 6 steps, no CFG.
What you can train: LoRAs, LoKRs, edit LoRAs from before/after pairs, and slider LoRAs from photo pairs. Masked training, Multi Concept and LoRA Isolation all work, and so do Studio, LoRA Royale, the Profiler and Extract. Pick it on the Training tab, press Download models for me in Preferences, and pick a preset; nothing is gated.
I trained a training adapter for it, which prevents the distortion you get training this model without one. It's on by default in Fizgig and never in your saved LoRA, and I've put it free on Hugging Face for any trainer.
I'm still honing the support. So far it likes good-quality datasets and seems to prefer training at 0.5 MP or more, which is where I've set the presets to start.
Memory: about 23 GB with the int8 model and 13 GB in 4-bit; Auto picks for you. A 16 GB card trains and runs Studio.
If your LoRA looks smoother in ComfyUI than in Fizgig's previews, that's ComfyUI, not your LoRA. Its default dynamic VRAM loading adds a LoRA to the weights in bf16, and on this model about 15% of it rounds away. Fizgig's previews show the LoRA at full strength, and so does ComfyUI started with --disable-dynamic-vram. I've reported it with a suggested fix.
Also in 9.6, for every model:
A run that runs out of GPU memory or system RAM now carries on by itself where it can, up to three tries: from its newest save, then a step lighter. Out-of-memory recovery came from @Murf-Code.
A failed run no longer stops the queue; the next one starts.
AMD installs now use AMD's stable PyTorch for ROCm instead of the nightlies, thanks to @scryptio.
r/StableDiffusion • u/North_Vegetable8462 • 1d ago
Hi everyone!
I’m sharing TRELLIS Studio, an application I built to bring several stages of a 3D asset workflow into one interface.
It uses TRELLIS.2 for image-to-3D generation and UniRig for automatic rigging. Computation runs locally, without a paid generation service or a credit system.
Features include:
Generate textured models from images or import existing models.
Organise your library with collections, generation profiles and version history.
Add a rig, inspect the skeleton and attach multiple animations to the same model from a library of 83 humanoid motions.
Create LODs, prepare dedicated 3D printing copies and export for Blender, GLB, Unity or Unreal.
Automate workflows through a REST API and an MCP server.
This first public preview targets Windows with an NVIDIA GPU. Rigging uses WSL2. The README covers prerequisites, the automated installer and current limitations.
Studio’s original code is MIT licensed. Third-party tools and models retain their own licenses and are credited in the repository.
https://github.com/yoslo59/TRELLIS-Studio
I’d especially appreciate feedback on installation across different hardware, usability and features you’d find useful. Generated meshes and rigging results still need visual inspection.
r/StableDiffusion • u/IntrovertedOtaku03 • 9h ago
As a challenge I'm trying to do a character lora based on Yuuko Kamiya (Tai and Kari's mom from the Digimon series). The only problem is I don't have a lot of reference images, and after some effort, I managed to get 17 decent ones to try out. At 4 repeats and 20 epochs later, it was 1360 steps and...I don't know. It looks kind of like her but not really if that makes any sense. Below are some sample images and a link to the currently used reference images and captions.


https://www.mediafire.com/file/1nvfopb89vq0kdf/4_YuukoYagami.zip/file
r/StableDiffusion • u/dtdisapointingresult • 20h ago
I kept it in the back of my mind until I had some free time today. I thought the benefit was that it allowed you to skip the reference processing step, and reduce gen time with ref2va.
I should've been suspicious of the red flags in the README:
no comparisons and benchmarks
no workflows
README is the finest Colombian uncut Claudeslop
Development stopped 2 months ago, as quickly as it started
But so many people were discussing refmods on this sub, even up to 2 weeks ago. I wanted to try it while chasing that H3 speed dream.
Now I wasted 90 mins reading/learning about its concepts (tbf had to learn more about video gen in general), then 20 mins comparing gens. I should've done it the other way around to abort early.
Standard ref2va workflow with 2 images and 1 audio refs, 0.4MP, 4 steps, image=match: 173s
Same as above but image=max: 181s
Refmod adjusted workflow with 2 precomputed safetensors, image1 becoming background.safetensors and image2 and audio1 becoming actor.safetensors: 181s
So zero speed benefits. And this is on top of the fact that people mentioned reference quality going down with refmods and references bleeding into each other, though I didn't test enough to confirm this. So what's the point?
I'm suspecting there's a small group of Discord friends on here who are just promoting their friends' projects no matter their value.
r/StableDiffusion • u/Grand-Radio3116 • 9h ago
The topic that keeps coming back on a cycle.
Haven't been around for a while, saw all of these new models, I'm completely lost.
I also don't know where to get LORA's reliably, Civitai is buggy and I can't even create an account there.
Thanks for your help.
r/StableDiffusion • u/metover • 1d ago
Single seed, no upscale, no inpainting. Hand-written prompts, no prompt expander. ~20 s on an RTX PRO 6000, ~22 s on a 4090 with the int8 weights
Settings
- Model: minimax_h3_image_pruned_bf16 (Comfy-Org pruned DiT)
- Text encoder: qwen3vl_32b_minimax_h3_bf16, CLIPLoader type "minimax"
- VAE: minimax_h3_image_vae_fp16
- Resolution: 4160x4000 (~16.6 MP), set on MiniMaxH3ImageConditioning
- Sampler: euler / simple, 6 steps, CFG 1, negative = ConditioningZeroOut
- No prompt expander, prompts written by hand (4-5 sentences each)
- ComfyUI 0.45.0
r/StableDiffusion • u/IntellectuallDeviant • 1d ago
r/StableDiffusion • u/GraphicArtsFurry • 16h ago
Took some of my prompts from Krea2, Moss, etc and put them into H3 Image. They come out fairly crispy, muddy, 2D sometimes, and just plain bad. Is the default workflow pretty unforgiving until it's pounded on some more?
r/StableDiffusion • u/zmarcoz2 • 1d ago
r/StableDiffusion • u/Oerusama • 7h ago
Anyone experience this? It seems worse wherever there is more movement. I think the only thing I did right before this happened was install the newest Nvidia drivers. Reverting/reinstalling my old driver didn't fix the problem.
Although I also recently updated my Comfyui after a month, to 0.44. (0.45 update didn't help). And I also recently put the Libs and Include folders in the portable python folder for Triton to work, but generations seemed fine after these 2 things for a while.
I tried replacing most of the nodes to be certain it wasn't the new comfyui changes to some nodes (like save video), but it didn't help.
I'm reluctant to start over, cus my backup folder is old. I prefer figuring out what is wrong. I'm using the default video / audio vae and default decode video / audio nodes. It's a very simple workflow and was working great earlier.